If the diffusion process is removing noise by predicting a final image and comparing it to the current one, why can't we just jump to the final predicted image? Or is the point that because its an iterative process, each noise step results in a different "final image" prediction?
To generate an image we start with a noise sample and take a step towards the _mean_ of the distribution of real images which would produce that noise sample when running the forward diffusion process. This step moves us towards the _mean_ of some distribution of real images and not towards a particular real image. But, as we take a bunch of small steps and gradually move back through the diffusion process, the effective distribution of real images over which this inverse diffusion prediction averages has lower and lower entropy, until it's effectively a specific real image, at which point we're done.
...but, the question is, why can't we take a big step and be at the end in one step.
Obviously a series of small steps gets you there, but the question was why you need to take small steps.
I feel like this is just a 'intuitive explanation' that doesn't actually do anything other than rephrase the question; "You take a series of small steps to reduce the noise in each step and end up with a picture with no noise".
The real reason is that big steps result in worse results (1); the model was specifically designed to be a series of small steps because when you take big steps, you end up with over fitting, where the model just generates a few outputs from any input.
Essentially, current approaches rely strongly on an assumption that the conditional we want to estimate in each step of the reverse diffusion process is approximately an isotropic Gaussian distribution. This assumption breaks down as you increase the size of the steps, and models which rely on the assumption also break down.
This is not directly related to overfitting. It is a fundamental aspect of how these models are designed and trained. If the architecture and loss function for training the inverse model were changed it would be possible to make an inverse model that inverts more steps of the forward diffusion process in a single go, but then the inverse model would need to become a full generative model on its own.
Hm. Why's that?
The only reason I mentioned over fitting is because that's literally what they say in the paper I linked, that the diffusion factor was selected to prevent over fitting.
...
I guess I don't really have a deep understanding of this stuff, but your explanation seems to be missing, specifically that noise is added to the latent each round, on a schedule (1), less noise each round.
that's what causes it to converge on a 'final' value; you're explicitly modifying the amount of additional noise you feed in. If you don't add any noise, you get nothing more from doing 1 step than you do from 10 or 50.
Right?
"as we take a bunch of small steps and gradually move back through the diffusion process, the effective distribution of real images over which this inverse diffusion prediction averages has lower and lower entropy"
I'm really not sure about that... :/
(1) - "For binomial diffusion, the discrete state space makes gradient ascent with frozen noise impossible. We instead choose the forward diffusion schedule β1···T to erase a constant fraction 1 T of the original signal per diffusion step, yielding a diffusion rate of βt = (T − t + 1)−1."
If t is small and the noise schedule for diffusion adds small noise at each step, then the inverse conditional p(x0 | xt) which we want to learn will be approximately a unimodal Gaussian. This is an intrinsic property of the forward diffusion process. When t is large, or the diffusion schedule adds a lot of noise at each step, the conditional p(x0 | xt) will be more complex and include a larger fraction of the images in the training set.
"If you don't add any noise, you get nothing more from doing 1 step than you do from 10 or 50." -- there are actually models which deterministically (approximately) integrate the reverse diffusion process SDE and don't involve any random sampling aside from the initial xT during generation.
For example, if t=T, where T is the total length of the diffusion process, then xt=xT is effectively an independent Gaussian sample and the inverse conditional p(x0 | xT) is simply p(x0) which is the distribution of the training data. In general, p(x0) is not a unimodal isotropic Gaussian. If it was, we could just model our training set by fitting the mean and (diagonal) covariance matrix of a Gaussian distribution.
"I'm really not sure about that... :/" -- the forward diffusion process initiated from x0 iteratively removes information about the starting point x0. Depending on the noise schedule, the rate at which information about x0 is removed by the addition of noise can vary. Whether we're in the continuous or discrete setting, this means the inverse conditional p(x0 | xt) will increase in entropy as t goes from 1 to T, where T is the max number of diffusion steps. So, when we generate an image by running the inverse diffusion process the conditional p(x0 | xt) will have shrinking entropy as t is now decreasing.
The quote (1) you reference is about how trying to directly optimize the noise schedule is more challenging when working with discrete inputs/latents. Whether the noise schedule is trained, as in their continuous case, or defined a priori, as in their discrete case, each step of forward diffusion removes information about the input and what I said about shrinking entropy of the conditional p(x0 | xt) as we run reverse diffusion holds. In the case of current SOTA diffusion models, I believe the noise schedules are set via hyperopt rather than optimized by SGD/ADAM/etc.
Because we're doing gradient descent. (No, seriously, it's turtles all the way down (or all the way up, considering we're at a higher level of abstraction here).)
We're trying to (quickly, in less than 100 steps) descend a gradient through a complex, irregular and heavily foggy 16384-dimensional landscape of smeared, distorted, and white-noise-covered images that kinda sorta look vaguely like what we want if you squint (well, if the neural network squints, anyway). If we try to take a big step, we don't descend the gradient faster; we fly off in a mostly random direction, clip through various proverbial cliffs, and probably end up somewhere higher up the gradient than we started.
You could also think of this as: We start with a terrible signal to noise ratio. So we need to average over very large areas to get any reasonable signal. But as we increase the signal, we can average over a smaller area to get the same signal-to-ratio.
In the beginning, we're averaging over large areas, so all the fine detail is lost. We just get 'might be a dog? maybe??'. What the network is doing is saying "if this a dog, there should be a head somewhere over here. So let me make it more like a head". Which improves the signal to noise ratio a bit.
After a few more steps, the signal is strong enough that we can get sufficient signal from smaller areas, so it starts saying 'head of a dog' in places. So the network will then start doing "Well, if this is a dog's head, there should be some eyes. Maybe two, but probably not three. And they'll be kinda somewhere around here".
Why do it this way?
Doing it this ways means the network doesn't need to learn "Here are all the ways dogs can look". Instead, it can learn a factored representation: A dog has a head and a body. The network only needs to learn a very fuzzy representation at this level. Then a head has some eyes and maybe a nose. Again, it only needs to learn a very fuzzy representation and (very) rough relative locations.
So it only when it get right down into fine detail that it actually needs to learn pixel perfect representation. But this is _way_ easier, because in small areas images have surprisingly very low entropy.
The 'text-to-image' bit is a just a twist on the basic idea. At the start when the network is going "dog? or it might be a horse?", we fiddle with the probabilities a bit so that the network starts out convinced there's a dog in there somewhere. At which point it starts making the most likely places look a little more like a dog.
There's a paper (which I can't remember the name of) that shows the process still works with different information removal operators, including one with a circle wipe, and one where it blends the original picture with a cat photo.
Also, this article describes CLIP being trained on text-image pairs, but Google's Imagen uses an off the shelf text model so that part doesn't seem to be needed either.
[1] https://arxiv.org/abs/2208.09392 [2] https://twitter.com/tomgoldsteincs/status/156250381442263040...
Typically, in generating the training data for diffusion models, we add noise incrementally to an image until it’s essentially all noise. Going backwards from almost all noise to the original images directly in one step is a pretty dubious proposition.
The diffusion process predicts the total noise that was added to the image. But that prediction isn't great and applying it immediately wouldn't result in a good output. So instead, the noise is multiplied by a small epsilon and then subtracted from the noisy image. That process is iterated to get to the final result.
1- Forward Diffusion (adding noise, and training the Unet to predict how much noise is added in each step)
2- Generating the image by denoising. This doesn't predict the final image, each step only predicts a small slice of noise (the removal of which leads to images similar to what the model encountered in step 1).
So it is indeed an iterative processes in that way, each step taking one step towards the final image.